stable-diffusion_skill

This skill helps you generate high-quality images from text prompts, perform image-to-image tasks, and optimize diffusion workflows with Stable Diffusion.
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Readme & install

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Installation

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npx veilstrat add skill orchestra-research/ai-research-skills --skill stable-diffusion

  • SKILL.md12.7 KB

Overview

This skill provides state-of-the-art text-to-image generation using Stable Diffusion models via the HuggingFace Diffusers library. It covers text-to-image, image-to-image, inpainting, ControlNet conditioning, LoRA adapters, and memory optimizations for GPU workflows. The content is practical and focused on getting reproducible, high-quality images with configurable pipelines and schedulers.

How this skill works

The skill orchestrates Stable Diffusion pipelines that convert text prompts into embeddings, run a denoising loop with a scheduler and a UNet/transformer noise predictor, and decode latents with a VAE into final images. It exposes utilities for image conditioning (image-to-image, masks, ControlNet) and adapter loading (LoRA) plus memory/precision settings to fit models on different hardware. You can swap schedulers, set guidance and negative prompts, and use torch generators for reproducible outputs.

When to use it

  • Generate photorealistic or stylized images from text prompts.
  • Transform or stylize existing images (image-to-image) or create variations.
  • Perform inpainting or outpainting to fill or extend image regions.
  • Apply spatial control using ControlNet (edges, pose, depth).
  • Prototype fast, low-step generation or produce high-quality outputs with SDXL.

Best practices

  • Use FP16/BF16 on compatible GPUs and enable model CPU offload when OOM occurs.
  • Tune guidance_scale (7–12 typical) and num_inference_steps for tradeoff between fidelity and speed.
  • Use negative_prompt to remove common artifacts and unwanted styles.
  • Swap to DPMSolverMultistep or UniPC for faster, high-quality convergence when reducing steps.
  • Enable attention/vae slicing and xFormers (if available) for large images or limited memory.

Example use cases

  • Create marketing visuals: generate multiple variations of a product scene for A/B testing.
  • Photo retouching: inpaint or replace objects in photos using masks.
  • Style transfer and character design: apply LoRA adapters to inject artistic styles.
  • Pose-to-image: use ControlNet openpose to generate consistent human poses.
  • Rapid prototyping: use LCM/LCM-scheduler and fused LoRA for 4–8 step preview renders.

FAQ

Set a torch.Generator with manual_seed and pass it via the generator parameter; keep seed, steps, scheduler and model weights constant.

What to do when I hit CUDA out-of-memory?

Enable model CPU offload, attention/vae slicing, use FP16, reduce batch size or image resolution, or run parts on CPU.

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stable-diffusion skill by orchestra-research/ai-research-skills | VeilStrat